# Text-to-Image Optimization

FLUX.1-dev model using DFloat11 lossless compression format, reducing VRAM usage by 30% while maintaining bit-level accuracy

LyCORIS adapter based on Stable Diffusion XL 1.0, focused on image generation tasks

A multimodal large language model fine-tuned based on Qwen2.5-VL-7B-Instruct, specifically optimized for text-to-image generation, capable of producing more detailed image descriptions

The GGUF quantized version of Lumina is a model specifically designed for generating high-quality images, supporting text-prompt-based generation of highly matched images.

An instruction-fine-tuned version based on the Pixtral-12B-2409 multimodal large language model, capable of generating richer detail descriptions for given images

An instruction-tuned version based on Qwen2-VL-7B-Instruct, focusing on generating more detailed image descriptions, optimized for text-to-image dataset creation.

Experimental GGUF-converted version of Flux.1-dev, featuring various mixed-precision quantization schemes

MaPO is a reference-free, energy-efficient, and memory-friendly alignment method for text-to-image diffusion models

This is a LoRA model based on Stable Diffusion XL, specifically designed for generating images with 'tears of joy' expressions.

This is a fine-tuned model based on Stable Diffusion XL, specializing in generating images with excellent lighting effects.

A Dreambooth model trained using TheLastBen's rapid DreamBooth technology, suitable for text-to-image tasks, employing stable diffusion.